# Art Style Transfer

This is a LoRA model that can convert ordinary images into metallic-style images with a black background and a 3D perspective.

LoRA fine-tuned model based on flux learning, specifically imitating Winslow Homer's artistic style

A training model based on flux learning, capable of generating oil paintings in the style of Paul Cezanne

ICEdit is an instructional image editing method based on large-scale diffusion transformers, achieving state-of-the-art editing effects with only 0.5% of training data and 1% of parameters.

This is a LoRA adaptation model based on Stable Diffusion XL, specifically designed for generating images in Van Gogh's artistic style.

A LoRA model trained on the Flex.1-alpha base model, specializing in generating magnificent natural landscape images, particularly excelling in sunset mountain scenes

This is a LoRA model for creating abstract art, developed based on the black-forest-labs/FLUX.1-dev foundation model.

IP-Adapter is an image prompt adapter for text-to-image diffusion models, capable of combining image prompts with text prompts to enhance the relevance and quality of generated images.

The IP Adapter is a text-compatible image prompt adapter for text-to-image diffusion models, capable of generating high-quality images through image prompts.

FLUX.1-dev is a text-to-image generation model focused on traffic-related image generation tasks.

An adapter model based on Stable Diffusion 1.5, focusing on image-to-image conversion, particularly excelling in maintaining image consistency under high-light conditions

KOALA is an advanced text-to-image generation model capable of producing high-quality images from textual descriptions.

This is a low-rank adaptation model based on Stable Diffusion XL, specifically designed to generate images in Studio Ghibli style.

An image editing model fine-tuned on Stable Diffusion XL (SDXL) using the InstructPix2Pix method, supporting image editing through natural language instructions.

DreamShaper is a powerful generative model focused on image generation and editing tasks.

Stable Diffusion Nano is a lightweight version fine-tuned on 128x128 images based on the stable diffusion model, suitable for rapid prototyping and experimentation.

ControlNet v1.1 is an extension model based on Stable Diffusion, adding extra conditions to control the image generation process.

A text-to-image generation model based on Stable Diffusion 1.5, focusing on artistic creation and self-expression themes

The first open-source Chinese-English bilingual Stable Diffusion model, trained on 20 million filtered Chinese image-text pairs

A Stable Diffusion model fine-tuned using Dreambooth technology, specifically designed to generate images in Hergé's style.